Stable Diffusion
Last updated
Last updated
Stable Diffusion is a text-to-image diffusion model capable of generating photo-realistic images given any text input. You can learn more about the model .
You can generate images using Stable Diffusion by using the "Create Prediction - Stable Diffusion" workflow action.
It requires the following inputs:
Prompt
Text
The text prompt you're using to create the image(s)
Height
Number (integer)
The height of the image(s) you want to create
Width
Number (integer)
The width of the image(s) you want to create
Number of Outputs
Number (integer)
The number of images you want to create with each prediction. Must be between 1 and 4.
Negative Prompt
Text
Things you do NOT want to see in the output image(s)
# Inference Steps
Number (integer)
The number of denoising steps used in the prediction (minimum: 1; maximum: 500). In general, the more steps you use the more detailed the output. However, the more steps you use the longer it will take to generate the prediction.
Guidance Scale
Number (decimal)
Scale for classifier-free guidance (minimum: 1; maximum: 20)
Seed
Number (integer)
Random seed. Leave blank to randomize the seed
Scheduler
Text (list of pre-defined options)
The type of scheduler you want to use.
Version
Text (list of pre-defined options)
The version of Stable Diffusion you want to use. I've found the default version (a9758cbfbd5f3c2094457d996681af52552901775aa2d6dd0b17fd15df959bef) works better than the newer version (ac732df83cea7fff18b8472768c88ad041fa750ff7682a21affe81863cbe77e4)